ckpt files) must be separately downloaded and are required to run Stable Diffusion. The Stable Diffusion 2. This video builds on the previous video which covered txt2img ( ) This video covers how to use Img2Img in Automat. Stable diffusion has been making huge waves recently in the AI and art communities (if you don’t know what that is feel free to check out this earlier post). 5 it/s. September 14, 2022 AI/ML. Predictions typically complete within 1 seconds. safetensors format. AIArtstable-diffusion-webuimodelsStable-diffusion768-v-ema. safetensors (5. . The program is tested to work on Python 3. I have searched the existing issues and checked the recent builds/commits What would your feature do ? with current technology would it be possible to ask the AI to generate a text from an image? in o. /webui. This extension adds a tab for CLIP Interrogator. A text-to-image generative AI model that creates beautiful images. Number of denoising steps. It. It’s a fun and creative way to give a unique twist to my images. Interrogation: Attempts to generate a list of words and confidence levels that describe an image. Stable Diffusion Hub. You can pull text from files, set up your own variables, process text through conditional functions, and so much more - it's like wildcards on steroids. Step 1: Go to DiffusionBee’s download page and download the installer for MacOS – Apple Silicon. Windows 11 Pro 64-bit (22H2) Our test PC for Stable Diffusion consisted of a Core i9-12900K, 32GB of DDR4-3600 memory, and a 2TB SSD. Generate high-resolution realistic images with AI. Stable diffusion image-to-text (SDIT) is an advanced image captioning model based on the GPT architecture and uses a diffusion-based training algorithm to improve stability and consistency during training. Type cmd. This model is a checkpoint merge, meaning it is a product of other models to create a product that derives. If you want to use a different name, use the --output flag. Discover stable diffusion Img2Img techniques & their applications. Introduction; Architecture; RequirementThe Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images. I created a reference page by using the prompt "a rabbit, by [artist]" with over 500+ artist names. Most people don't manually caption images when they're creating training sets. The CLIP Interrogator is a prompt engineering tool that combines OpenAI's CLIP and Salesforce's BLIP to optimize text prompts to match a given image. Is there an alternative. 学習元のモデルが決まったら、そのモデルを使った正則化画像を用意します。 ここも必ず必要な手順ではないので、飛ばしても問題ありません。Stable Diffusion. Stable Diffusion img2img support comes to Photoshop. r/StableDiffusion. StableDiffusion. A buddy of mine told me about it being able to be locally installed on a machine. 1 (diffusion, upscaling and inpainting checkpoints) 🆕 Now available as a Stable Diffusion Web UI Extension! 🆕. Replicate makes it easy to run machine learning models in the cloud from your own code. Stable DiffusionはNovelAIやMidjourneyとはどう違うの？ Stable Diffusionを簡単に使えるツールは結局どれを使えばいいの？ 画像生成用のグラフィックボードを買うならどれがオススメ？ モデルのckptとsafetensorsって何が違うの？ モデルのfp16・fp32・prunedって何？本教程需要一些AI绘画基础，并不是面对0基础人员，如果你没有学习过stable diffusion的基本操作或者对Controlnet插件毫无了解，可以先看看秋葉aaaki等up的教程，做到会存放大模型，会安装插件并且有基本的视频剪辑能力。-----一、准备工作This issue is a workaround for a security vulnerability. 1M runs. Stable Diffusion web UIをインストールして使えるようにしておく。 Stable Diffusion web UI用のControlNet拡張機能もインストールしておく。 この2つについては下記の記事でやり方等を丁寧にご説明していますので、まだ準備ができていないよという方はそちらも併せて. As of June 2023, Midjourney also gained inpainting and outpainting via the Zoom Out button. Type and ye shall receive. The weights were ported from the original implementation. The CLIP interrogator has two parts: one is the BLIP model, which takes on the function of decoding and reasoning about the text description. GitHub. Share generated images with LAION for improving their dataset. 1 I use this = oversaturated, ugly, 3d, render, cartoon, grain, low-res, kitsch, black and white. この記事では と呼ばれる手法で、画像からテキスト（プロンプト）を取得する方法を紹介します。. . r/sdnsfw Lounge. #. AI画像生成士. ckpt for using v1. this Stable diffusion model i have fine tuned on 1000 raw logo png/jpg images of of size 128x128 with augmentation. 1 images, the RTX 4070 still plugs along at over nine images per minute (59% slower than 512x512), but for now AMD's fastest GPUs drop to around a third of. safetensors files from their subfolders if they’re available in the model repository. On SD 2. Image to text, img to txt. Fix it to look like the original. 打开stable-diffusion-webuimodelsstable-diffusion目录，此处为各种模型的存放处。 需要预先存放一个模型才能正常使用。 3. More posts you may like r/selfhosted Join • 13. Download: Installation: Extract anywhere (not a protected folder - NOT Program Files - preferrably a short custom path like D:/Apps/AI/), run StableDiffusionGui. img2txt. Full model fine-tuning of Stable Diffusion used to be slow and difficult, and that's part of the reason why lighter-weight methods such as Dreambooth or Textual Inversion have become so popular. Stable Diffusion is a concealed text-to-image diffusion model, capable of generating photorealistic images from any textual input, fosters independent flexibility in producing remarkable visuals. 手順3：学習を行う. 0 model. 0 was released in November 2022 and has been entirely funded and developed by Stability AI. But the width, height and other defaults need changing. September 14, 2022 AI/ML. To quickly summarize: Stable Diffusion (Latent Diffusion Model) conducts the diffusion process in the latent space, and thus it is much faster than a pure diffusion model. img2txt online. 前提：Stable. 零基础学会Stable Diffusion，这绝对是你看过的最容易上手的AI绘画教程 | SD WebUI 保姆级攻略，一站式入门AI绘画！Midjourney胎教级入门指南！普通人也能成为设计师，图片描述的答题技巧，Stable Diffusion 反推提示词的介绍及运用（cilp、deepbooru） 全流程教程（教程合集. StabilityAI’s Stable Video Diffusion (SVD), image to video Updated 4 hours ago 173 runs sdxl A text-to-image generative AI model that creates beautiful images Updated 2 weeks, 2 days ago 20. The idea behind the model was derived from my ReV Mix model. It may help to use the inpainting model, but not. Img2Txt. See the complete guide for prompt building for a tutorial. So the style can match the original. However, at the time he installed it only one . Introduction. jkcarney commented Jun 30, 2023. In this video we'll walk through how to run Stable Diffusion img2img and txt2img using AMD GPU on Windows operating system. r/StableDiffusion •. A checker for NSFW images. Wait a few moments, and you'll have four AI-generated options to choose from. This model runs on Nvidia T4 GPU hardware. 以下方式部署的stable diffusion ui仅会使用CPU进行计算，在没有gpu加速的情况下，ai绘图会占用 非常高（几乎全部）的CPU资源 ，并且绘制单张图片的 时间会比较长 ，仅建议CPU性能足够强的情况下使用（作为对比参考，我的使用环境为笔记本平台的5900HX，在默认参数. Run time and cost. Then we design a subject representation learning task, called prompted. Join. You'll have a much easier time if you generate the base image in SD, add in text with a conventional image editing program. . Easy Prompt SelectorのYAMLファイルは「stable-diffusion-webuiextensionssdweb-easy-prompt-selector ags」の中にあります。 「. 手順2：「gui. Improving image generation at different aspect ratios using conditional masking during training. ai says it can double the resolution of a typical 512×512 pixel image in half a second. 160 upvotes · 39 comments. The domain img2txt. It is common to use negative embeddings for anime. So once you find a relevant image, you can click on it to see the prompt. Predictions typically complete within 27 seconds. Dreambooth is considered more powerful because it fine-tunes the weight of the whole model. py", line 222, in load_checkpoint raise RuntimeError('checkpoint url or path is invalid') The text was updated successfully, but these errors were encountered: All reactions. Find your API token in your account settings. Also, because the Payload source code is fully written in. 08:41. You can receive up to four options per prompt. It includes every name I could find in prompt guides, lists of. with current technology would it be possible to ask the AI to generate a text from an image? in order to know what technology could describe the image, a tool for AI to describe the image for us. With its 860M UNet and 123M text encoder. Qualcomm has demoed AI image generator Stable Diffusion running locally on a mobile in under 15 seconds. I. Další příspěvky na téma Stable Diffusion. Inside your subject folder, create yet another subfolder and call it output. BLIP-2 is a zero-shot visual-language model that can be used for multiple image-to-text tasks with image and image and text prompts. 上个月做了安卓和苹果手机用远端sd进行跑图的几个demo，整体流程很简单. The Payload Config. By default this will display the “Stable Diffusion Checkpoint” drop down box which can be used to select the different models which you have saved in the “stable-diffusion-webuimodelsStable-diffusion” directory. 画像→テキスト（img2txt）は、Stable Diffusionにも採用されている CLIP という技術を使います。 CLIPは簡単にいうと、単語をベクトル化（数値化）することで計算できるように、さらには他の単語と比較できるようにするものです。Run time and cost. 🖊️ sd-2. Hey there! I’ve been doing some extensive tests between diffuser’s stable diffusion and AUTOMATIC1111’s and NMKD-SD-GUI implementations (which both wrap the CompVis/stable-diffusion repo). 10. Mage Space has very limited free features, so it may as well be a paid app. この記事ではStable diffusionが提供するAPIを経由して、. ChatGPT is aware of the history of your current conversation. Using the above metrics helps evaluate models that are class-conditioned. Pak jsem si řekl, že zkusím img2txt a ten vytvořil. 89 GB) Safetensors Download ProtoGen x3. This is a builtin feature in webui. You can use 6-8 GB too. Yodayo gives you more free use, and is 100% anime oriented. You should see the message. 10. Installing. For example, DiT. Just two. Starting from a random noise, the picture is enhanced several times and the final result is supposed to be as close as possible to the keywords. Image-to-Text Transformers. What is Img2Img in Stable Diffusion Setting up The Software for Stable Diffusion Img2img How to Use img2img in Stable Diffusion Step 1: Set the. The latest stability ai release is 2. 1:7860" or "localhost:7860" into the address bar, and hit Enter. creates original designs within seconds. When it comes to speed to output a single image, the most powerful. This model runs on Nvidia A100 (40GB) GPU hardware. Then you can either mask the face and choose inpaint unmasked, or select only the parts you want changed and inpaint masked. Generated in -4480634. 9 conda activate 522-project # install torch 2. Start with installation & basics, then explore advanced techniques to become an expert. But in addition, there's also a Negative Prompt box where you can preempt Stable Diffusion to leave things out. Check it out: Stable Diffusion Photoshop Plugin (0. エイプリルフールのネタとして自分の長年使ってきたTwitterアイコンを変えるのを思いついたはいいものの、素材をどうするかということで流行りのStable Diffusionでつくってみました。. Change the sampling steps to 50. (Optimized for stable-diffusion (clip ViT-L/14)) Public; 2. ai and more. openai. com uses a Commercial suffix and it's server(s) are located in N/A with the IP number 104. I am still new to Stable Diffusion, but I still managed to get an art piece with text, nonetheless. Scroll to the bottom of the notebook to the Prompts section near the very bottom of the notebook. By default, Colab notebooks rely on the original Stable Diffusion which comes with NSFW filters. stable diffusion webui 脚本使用方法（下），人脸编辑还不错. The following resources can be helpful if you're looking for more. The result can be viewed on 3D or holographic devices like VR headsets or lookingglass display, used in Render- or Game- Engines on a plane with a displacement modifier, and maybe even 3D printed. 5 base model. Height. 0 和 2. Text to image generation. Para ello vam. Whilst the then popular Waifu Diffusion was trained on SD + 300k anime images, NAI was trained on millions. Stability AI’s Stable Diffusion, high fidelity but capable of being run on off-the-shelf consumer hardware, is now in use by art generator services like Artbreeder, Pixelz. Roboti na kole. Its installation process is no different from any other app. You can use them to remove specific elements, styles, or. js client: npm install replicate. Contents. After applying stable diffusion techniques with img2img, it's important to. 本文接下来就会从效果及原理两个部分介绍Diffusion Model，具体章节如下：. /. Check it out: Stable Diffusion Photoshop Plugin (0. Those are the absolute minimum system requirements for Stable Diffusion. Model card Files Files and versions Community Train. 2. stable diffusion webui 脚本使用方法（上）. You can also upload and replicate non-AI generated images. Stable diffusion is an open-source technology. • 5 mo. For 2. It came out gibberish though. Using stable diffusion and these prompts hand-in-hand, you can easily create stunning and high-quality logos in seconds without needing any design experience. (with < 300 lines of codes!) (Open in Colab) Build. photo of perfect green apple with stem, water droplets, dramatic lighting. bat (Windows Batch File) to start. Waifu Diffusion 1. Stable Horde for Web UI. The script outputs an image file based on the model's interpretation of the prompt. To use img2txt stable diffusion, all you need to do is provide the path or URL of the image you want to convert. 6 The Stable Diffusion 2 repository implemented all the servers in gradio and streamlit model-type is the type of image modification demo to launch For example, to launch the streamlit version of the image upscaler on the model created in the original step (assuming the x4-upscaler-ema. Functioning as image viewers for the terminal, chafa and catimg have only been an integral part of a stable release of the Linux distribution since Debian GNU/Linux 10. (You can also experiment with other models. Let's dive in deep and learn how to generate beautiful AI Art based on prom. . This may take a few minutes. Search by model Stable Diffusion Midjourney ChatGPT as seen in. Also there is post tagged here where all the links to all resources are. Given a (potentially crude) image and the right text prompt, latent diffusion. The goal of this article is to get you up to speed on stable diffusion. The most popular image-to-image models are Stable Diffusion v1. 0 的过程，包括下载必要的模型以及如何将它们安装到. Step 3: Clone web-ui. (Optimized for stable-diffusion (clip ViT-L/14)) Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from CompVis, Stability AI and LAION. Jolly-Theme-7570. 解析度拉越高，所需算圖時間越久，VRAM 也需要更多、甚至會爆顯存，因此提高的解析度有上限. This process is called "reverse diffusion," based on math inspired. For those of you who don’t know, negative prompts are things you want the image generator to exclude from your image creations. ago. Create multiple variants of an image with Stable Diffusion. . 4. This controls the resolution which an image is initially generated at. These encoders are trained to maximize the similarity of (image, text) pairs via a contrastive loss. Uncrop. The default we use is 25 steps which should be enough for generating any kind of image. SD教程•重磅更新！. I wanted to report some observations and wondered if the community might be able to shed some light on the findings. Creating venv in directory C:UsersGOWTHAMDocumentsSDmodelstable-diffusion-webuivenv using python "C:UsersGOWTHAMAppDataLocalProgramsPythonPython310python. Are there options for img2txt and txt2txt I'm working on getting GPT-J and stable diffusion working on proxmox and it's just amazing, now I'm wondering what else can this tech do ? And by txt2img I would expect you feed our an image and it tells you in text what it sees and where. Documentation is lacking. Check out the Quick Start Guide if you are new to Stable Diffusion. Roughly: Use IMG2txt. 0) Watch on. Predictions typically complete within 2 seconds. All the training scripts for text-to-image finetuning used in this guide can be found in this repository if you’re interested in taking a closer look. Next, copy your API token and authenticate by setting it as an environment variable: export REPLICATE_API_TOKEN=<paste-your-token-here>. Set the batch size to 4 so that you can. 12GB or more install space. It generates accurate, diverse and creative captions for images. 0-base. r/StableDiffusion. If you don't like the results, you can generate new designs an infinite number of times until you find a logo you absolutely love! Watch It In Action. The results from the Stable Diffusion and Kandinsky models vary due to their architecture differences and training process; you can generally expect SDXL to produce higher quality images than Stable Diffusion v1. 2. The learned concepts can be used to better control the images generated from text-to-image. Transform your doodles into real images in seconds. Generate and Run Olive Optimized Stable Diffusion Models with Automatic1111 WebUI on AMD GPUs. img2txt stable diffusion. Stable Doodle. 152. ckpt for using v1. Create beautiful images with our AI Image Generator (Text to Image) for free. It can be done because I saw it with. The generated image will be named img2img-out. stable-diffusion-img2img. Hiresは「High Resolution」の略称で高解像度という意味を持ち、fixは「修正・変更」を指します。. photo of perfect green apple with stem, water droplets, dramatic lighting. ago. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Enjoy . If there is a text-to-image model that can come very close to Midjourney, then it’s Stable Diffusion. 😉. Goodbye Babel, generated by Andrew Zhu using Diffusers in pure Python. On SD 2. . 002. There is no rule here - the more area of the original image is covered, the better match. Here is how to generate Microsoft Olive optimized stable diffusion model and run it using Automatic1111 WebUI: Open Anaconda/Miniconda Terminal. Caption. As with all things Stable Diffusion, the checkpoint model you use will have the biggest impact on your results. Stable Diffusion Prompts Generator helps you. London- and California-based startup Stability AI has released Stable Diffusion, an image-generating AI that can produce high-quality images that look as if they were. ckpt checkpoint was downloaded), run the following: Technical details regarding Stable Diffusion samplers, confirmed by Katherine: - DDIM and PLMS are originally the Latent Diffusion repo DDIM was implemented by CompVis group and was default (slightly different update rule than the samplers below, eqn 15 in DDIM paper is the update rule vs solving eqn 14's ODE directly) While Stable Diffusion doesn't have a native Image-Variation task, the authors recreated the effects of their Image-Variation script using the Stable Diffusion v1-4 checkpoint. Sort of new here. All stylized images in this section is generated from the original image below with zero examples. We follow the original repository and provide basic inference scripts to sample from the models. By simply replacing all instances linking to the original script with the script that has no safety filters, you can easily achieve generate NSFW images. This is a repo providing same stable diffusion experiments, regarding textual inversion task and captioning task pytorch clip captioning-images img2txt caption-generation caption-generator huggingface latent-diffusion stable-diffusion huggingface-diffusers latent-diffusion-models textual-inversion VGG16 Guided Stable Diffusion. A negative prompt is a way to use Stable Diffusion in a way that allows the user to specify what he doesn’t want to see, without any extra input. ; Mind you, the file is over 8GB so while you wait for the download. To try it out, tune the H and W arguments (which will be integer-divided by 8 in order to calculate the corresponding latent size), e. For the rest of this guide, we'll either use the generic Stable Diffusion v1. You are welcome to try our free online Stable Diffusion based image generator at It supports img2img generation, including sketching of the initial image :) Cool site. The train_text_to_image. Want to see examples of what you can build with Replicate? Check out our showcase. 0. langchain load local huggingface model example in python The following describes an example where a rough sketch. This model runs on Nvidia T4 GPU hardware. A surrealist painting of a cat by Salvador Dali/r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. 使用MediaPipe的面部网格注释器的修改输出，在LAION-Face数据集的一个子集上训练了ControlNet，以便在生成面部图像时提供新级别的控. Useful resource. Linux: run the command webui-user. 4); stable_diffusion (v1. Use your browser to go to the Stable Diffusion Online site and click the button that says Get started for free. First, your text prompt gets projected into a latent vector space by the. Public. Stable Diffusion - Image to Prompts Run 934. File "C:UsersGros2stable-diffusion-webuildmmodelslip. Stable Diffusion v1. 生成按钮下有一个 Interrogate CLIP，点击后会下载 CLIP，用于推理当前图片框内图片的 Prompt 并填充到提示词。 CLIP 询问器有两个部分：一个是 BLIP 模型，它承担解码的功能，从图片中推理文本描述。 The Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images. 1 I use this = oversaturated, ugly, 3d, render, cartoon, grain, low-res, kitsch, black and white. img2txt ascii. • 7 mo. 21. Works in the same way as LoRA except for sharing weights for some layers. Aspect ratio is kept but a little data on the left and right is lost. Type a question in the input box at the bottom to start a conversation. portrait of a beautiful death queen in a beautiful mansion painting by craig mullins and leyendecker, studio ghibli fantasy close - up shot. Stable Diffusion XL. The original implementation had two variants: one using a ResNet image encoder and the other. Steps. 2. 04 for arm 32 bitIt's wild to think Photoshop has a Stable Diffusion Text to A. You'll have a much easier time if you generate the base image in SD, add in text with a conventional image editing program. Prompt string along with the model and seed number. env. 20. World of Warcraft? Návrat ke kostce, a vyšel neuvěřitelně. Second day with Animatediff, SD1. One of the most amazing features is the ability to condition image generation from an existing image or sketch. sh in terminal to start. Stable Diffusion一键AI绘画、捏脸改图换背景，从安装到使用. Available values: 21, 31, 41, 51. GitHub. CLIP Interrogator extension for Stable Diffusion WebUI. So the Unstable Diffusion. plugin already! NOTE: Once installed, you will be able to generate images without a subscrip. If you don't like the results, you can generate new designs an infinite number of times until you find a logo you absolutely love! Watch It In Action. . It is a parameter that tells the Stable Diffusion model what not to include in the generated image. Install the Node. See the SDXL guide for an alternative setup with SD. Embeddings (aka textual inversion) are specially trained keywords to enhance images generated using Stable Diffusion. 31 votes, 370 comments. 667 messages. Get prompts from stable diffusion generated images. This example was created by a different version, rmokady/clip_prefix_caption:d703881e. Checkpoints (. There have been a few recent threads about approaches for this sort of thing and I'm always interested to see what new ideas people have. generating img2txt with the new v2. Appendix A: Stable Diffusion Prompt Guide. they converted to a. LoRAを使った学習のやり方. Set image width and height to 512. The inspiration was simply the lack of any Emiru model of any sort here. This endpoint generates and returns an image from a text passed in the request body. Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. Overview Stable Diffusion V3 APIs Text2Image API generates an image from a text prompt. hatenablog. 0. ago. Hi, yes you can mix two even more images with stable diffusion. ago. To run the same text-to-image prompt as in the notebook example as an inference job, use the following command: trainml job create inference "Stable Diffusion. Get inspired with Kiwi Prompt's stable diffusion prompts for clothes. It can be used in combination with. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Authors: Christoph Schuhmann, Richard Vencu, Romain Beaumont, Theo Coombes, Cade Gordon, Aarush Katta, Robert Kaczmarczyk, Jenia JitsevFirst, choose a diffusion model on promptoMANIA and put down your prompt or the subject of your image. ago. A diffusion model, which repeatedly "denoises" a 64x64 latent image patch. 16:17. Below are some of the key features: – User-friendly interface, easy to use right in the browser – Supports various image generation options like size, amount, mode,. r/StableDiffusion •. 使用管理员权限打开下图应用程序. Stable Diffusion without UI or tricks (only take off filter xD). Text prompt with description of the things you want in the image to be generated. 因為是透過 Stable Diffusion Model 算圖，除了放大解析度外，還能增加細部細節！. CLIP Interrogator extension for Stable Diffusion WebUI. By my understanding, a lower value will be more "creative" whereas a higher value will adhere more to the prompt. 5 it/s (The default software) tensorRT: 8 it/s. Software to use SDXL model. However, there’s a twist. This parameter controls the number of these denoising steps. • 1 yr. 5. 以 google. Just go to this address and you will see and learn: Fine-tune Your AI Images With These Simple Prompting Techniques - Stable Diffusion Art (stable-diffusion-art. Syntax: cv2. In the dropdown menu, select the VAE file you want to use. This model runs on Nvidia A40 (Large) GPU hardware. Hosted on Banana 🍌. Resize and fill: This will add in new noise to pad your image to 512x512, then scale to 1024x1024, with the expectation that img2img will. Request --request POST '\ Run time and cost. A snaha vytvořit obrázek…Anime embeddings. py", line 222, in load_checkpoint raise RuntimeError('checkpoint url or path is invalid') The text was updated successfully, but these errors were encountered: Most common negative prompts according to SD community. The client will automatically download the dependency and the required model. But the […]Stable Diffusion是2022年發布的深度學習 文本到图像生成模型。 它主要用於根據文本的描述產生詳細圖像，儘管它也可以應用於其他任務，如內補繪製、外補繪製，以及在提示詞指導下產生圖生圖的转变。. Get the result. 0. Img2Prompt. josemuanespinto. Notice there are cases where the output is barely recognizable as a rabbit. Diffusers dreambooth runs fine with --gradent_checkpointing and adam8bit, 0. 2. methexis-inc / img2prompt. stability-ai. Stable Diffusion consists of three parts: A text encoder, which turns your prompt into a latent vector.